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my-korean-stable-diffusion-v1-5

stable-diffusion-x4-upscaler

This model card focuses on the model associated with the Stable Diffusion Upscaler, available here. This model is trained for 1.25M steps on a 10M subset of LAION containing images `>2048x2048`. The model was trained on crops of size `512x512` and is a text-guided latent upscaling diffusion model. In addition to the textual input, it receives a `noise_level` as an input parameter, which can be used to add noise to the low-resolution input according to a predefined diffusion schedule. @InProceedings{Rombach_2022_CVPR, author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn}, title = {High-Resolution Image Synthesis With Latent Diffusion Models}, booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)}, month = {June}, year = {2022}, pages = {10684-10695} } **Notes**: - Despite not being a dependency, we highly recommend you to install xformers for memory efficient attention (better performance) - If you have low GPU RAM available, make sure to add a `pipe.enable_attention_slicing()` after sending it to `cuda` for less VRAM usage (to the cost of speed)

juggernaut-reborn

Coding in PHP/Node/Java etc? Have a look at docs for more code examples: View docs Try model for free: Generate Images payload = json.dumps({ "key": "your_api_key", "model_id": "juggernaut-reborn", "prompt": "ultra realistic close up portrait ((beautiful pale cyberpunk female with heavy black eyeliner)), blue eyes, shaved side haircut, hyper detail, cinematic lighting, magic neon, dark red city, Canon EOS R3, nikon, f/1.4, ISO 200, 1/160s, 8K, RAW, unedited, symmetrical balance, in-frame, 8K", "negative_prompt": "painting, extra fingers, mutated hands, poorly drawn hands, poorly drawn face, deformed, ugly, blurry, bad anatomy, bad proportions, extra limbs, cloned face, skinny, glitchy, double torso, extra arms, extra hands, mangled fingers, missing lips, ugly face, distorted face, extra legs, anime", "width": "512", "height": "512", "samples": "1", "num_inference_steps": "30", "safety_checker": "no", "enhance_prompt": "yes", "seed": None, "guidance_scale": 7.5, "multi_lingual": "no", "panorama": "no", "self_attention": "no", "upscale": "no", "embeddings": "embeddings_model_id", "lora": "lora_model_id", "webhook": None, "track_id": None })

controlnet-canny-sdxl-1.0

Hello, I am very happy to announce the controlnet-canny-sdxl-1.0 model, **a very powerful controlnet that can generate high resolution images visually comparable with midjourney**. The model was trained with large amount of high quality data(over 10000000 images), with carefully filtered and captioned(powerful vllm model). Besides, useful tricks are applied during the training, including date augmentation, mutiple loss and multi resolution. With only 1 stage training, the performance outperforms the other opensource canny models ([diffusers/controlnet-canny-sdxl-1.0], [TheMistoAI/MistoLine]). I release it and hope to advance the application of stable diffusion models. Canny is one of the most important ControlNet series models and can be applied to many jobs associated with drawing and designing. - **Developed by:** xinsir - **Model type:** ControlNet_SDXL - **License:** apache-2.0 - **Finetuned from model [optional]:** stabilityai/stable-diffusion-xl-base-1.0 prompt: A closeup of two day of the dead models, looking to the side, large flowered headdress, full dia de Los muertoe make up, lush red lips, butterflies, flowers, pastel colors, looking to the side, jungle, birds, color harmony , extremely detailed, intricate, ornate, motion, stunning, beautiful, unique, soft lighting

Sakura-MiniMax-H3-GGUF

Community quantizations of the pruned FL2VA denoiser (the diffusion transformer, about 20 B parameters) of MiniMax-H3, the model that generates video with a synchronized stereo audio track. Made with stable-diffusion.cpp. Independent work, not made or endorsed by MiniMax or Comfy-Org. This repository holds the **denoiser** and unmodified copies of the two **VAEs**. The Qwen3-VL text encoder is a separate file: an abliterated ("uncensored") build is available as a separate repository, Sakura-MiniMax-H3-Uncensored-GGUF; unchanged encoders in GGUF exist from other community members (for example leejet/MiniMax-H3-GGUF). ComfyUI: not tested by us; the community GGUF loader fork named in leejet's repository is the one to try.

stable-diffusion-v1-5-GGUF

**Model creator**: Runway ML **Original model**: stable-diffusion-v1-5 **GGUF quantization**: based on stable-diffusion.cpp ac54e that patched by llama-box. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. For more information about how Stable Diffusion functions, please have a look at 🤗's Stable Diffusion blog. The **Stable-Diffusion-v1-5** checkpoint was initialized with the weights of the Stable-Diffusion-v1-2 checkpoint and subsequently fine-tuned on 595k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling.

stable-diffusion-inpainting-gguf

HunyuanWorld-1

We would like to thank the contributors to the Stable Diffusion, FLUX, diffusers, HuggingFace, Real-ESRGAN, ZIM, GroundingDINO, MoGe, Worldsheet, WorldGen repositories, for their open research.

stable-diffusion-2-base

minecraft-skin-generator-sdxl

Monadical is pleased to announce the official release of the Minecraft Skin Generator XL model. We had previously released the Minecraft Skin Generator model based upon Stable Diffusion 2. This new model offers significant improvements over the last generation of models. 1. **Upgrade to Stable Diffusion XL** - Our model is now based upon the Stable Diffusion XL model, which greatly improves the quality of generated skins when compared to previous models. 1. **Transparent Layer Support** - The new model now supports the transparency layer in the hair and helmet section of the skin. * 'Saul Goodman from the TV show "Better Call Saul"'