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Flux

stable-diffusion-xl-base-1.0-GGUF

**Model creator**: Stability AI **Original model**: stable-diffusion-xl-base-1.0 **GGUF quantization**: based on stable-diffusion.cpp ac54e that patched by llama-box. **VAE From**: madebyollin/sdxl-vae-fp16-fix. SDXL consists of an ensemble of experts pipeline for latent diffusion: In a first step, the base model is used to generate (noisy) latents, which are then further processed with a refinement model (available here: https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0/) specialized for the final denoising steps. Note that the base model can be used as a standalone module. Alternatively, we can use a two-stage pipeline as follows: First, the base model is used to generate latents of the desired output size. In the second step, we use a specialized high-resolution model and apply a technique called SDEdit (https://arxiv.org/abs/2108.01073, also known as "img2img") to the latents generated in the first step, using the same prompt. This technique is slightly slower than the first one, as it requires more function evaluations.

stable-diffusion-3.5-large-DF11

stable-diffusion-3.5-medium-GGUF

control_v11p_sd15_mlsd

**Controlnet v1.1** is the successor model of Controlnet v1.0 and was released in lllyasviel/ControlNet-v1-1 by Lvmin Zhang. This checkpoint is a conversion of the original checkpoint into `diffusers` format. It can be used in combination with **Stable Diffusion**, such as runwayml/stable-diffusion-v1-5. @misc{zhang2023adding, title={Adding Conditional Control to Text-to-Image Diffusion Models}, author={Lvmin Zhang and Maneesh Agrawala}, year={2023}, eprint={2302.05543}, archivePrefix={arXiv}, primaryClass={cs.CV} }

FLUX.1-schnell-gguf

The license of the quantized files follows the license of the original model: - FLUX.1-schnell: apache-2.0 Run FLUX using stable-diffusion.cpp with a GPU that has 6GB or even 4GB of VRAM: https://github.com/leejet/stable-diffusion.cpp/blob/master/docs/flux.md

stable-diffusion-1.5

This is an archival re-upload of Stable Diffusion v1.5, originally at https://huggingface.co/runwayml/stable-diffusion-v1-5 until RunwayML took down that page. This model is from 2022, and is several major generational upgrades behind, it is being preserved here for technical & accessibility reasons (eg legacy model testing). https://huggingface.co/Comfy-Org/stable-diffusion-v1-5-archive/blob/main/v1-5-pruned-emaonly.safetensors is the exact original hash-identical model as uploaded by RunwayML.

stable-diffusion-3-medium-tensorrt

tiny-random-stable-diffusion-xl-refiner

zero123-xl-diffusers

_Note: This section is originally taken from the Stable Diffusion v2 model card, but applies in the same way to Zero-1-to-3._ The model is intended for research purposes only. Possible research areas and tasks include: - Safe deployment of large-scale models. - Probing and understanding the limitations and biases of generative models. - Generation of artworks and use in design and other artistic processes. - Applications in educational or creative tools. - Research on generative models.